5 model. After an entire weekend reviewing the material, I think (I hope!) I got the implementation right: As the title says, I included ControlNet XL OpenPose and FaceDefiner models. Basically a bunch of junk so that I can perfect the image. The refiner should definitely NOT be used as the starting point model for text2img. Misconfiguring nodes can lead to erroneous conclusions, and it's essential to understand the correct settings for a fair assessment. People using utilities like Textual Inversion and DreamBooth have been able to solve the problem in narrow use cases, but to the best of my knowledge there isn't yet a reliable solution to make on-model characters without just straight up hand-holding the AI. This accuracy allows much more to be done to get the perfect image directly from text, even before using the more advanced features or fine-tuning that Stable Diffusion is famous for. 2), (isometric 3d art of floating rock citadel:1), cobblestone, flowers, verdant, stone, moss, fish pool, (waterfall:1. 0-RC. 0, VAE hash example of workflow: Prompt: full body photo of beautiful age 18 girl, elf ears, blonde hair, freckles, sexy, beautiful, BREAK hiding behind a tree in the forest /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. This one feels like it starts to have problems before the effect can Is there an explanation for how to use the refiner in ComfyUI? You can just use someone elses workflow of 0. It's called Family pack, get with the times old man! Dude is corn. I've tried this article, but the result does not give me what I wanted. It'll be at the top though, not where it used to be. A TON of budget of commercial shoots is location-based. I guess an important thing for the quality (when you mention without finetunes) is that this time, the base model is finetuned by Lykon, the number 1 model creator on civitai. But it's reasonably clean to be used as a learning tool, which is and will It's amazing - I can get 1024x1024 SDXL images in ~40 seconds at 40 iterations euler A with base/refiner with the medvram-sdxl flag enabled now. img2img API with inpainting. They also have an SDXL Lora that kinda adds some contrast. 0 Refine. I also automated the split of the diffusion steps between the Base and the I thought my gaming would be at least a lot better than my 2070 super 8gb. 9 vae in there. Kohya Deepshrink is based on Scalecrafter research. These comparisons are useless without knowing your workflow. I had to use clip interrogator on Replicate because it gives me errors when using it locally. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. For example, if you wanted a great image of a person in a firefighter outfit, you could add a specific extra ‘embedding’ model trained on images of firefighter outfits. On the other hand tin my experience the SD3 renders doesn't mix very well with refiners so what you obtain is almost a dead end. I was really like stable Cascade mixed with a 1. This brings back memories of the first time that I use Stable Diffusion myself. Can take a while, on average I need two or the dalle2 inpainting prompts to get them fixed. 3~0. Thanks tons! Accidentally used the refiner model to generate images. 0 for ComfyUI - Now with Face Swapper, Prompt Enricher (via OpenAI), Image2Image (single images and batches), FreeU v2, XY Plot, ControlNet and ControlLoRAs, SDXL Base + Refiner, Hand Detailer, Face Detailer, Upscalers, ReVision, etc. Create a Load Checkpoint node, in that node select the sd_xl_refiner_0. The first is PixArt Sigma with no refinement and the second is after a . 519K subscribers in the StableDiffusion Make sure you have: Settings -> Stable Diffusion -- > "Maximum number of checkpoints loaded at the same time" set to 2 so it wont unload and reload the model for each pass. I will first try out the newest sd. 0. Inpainting is almost always needed to fix the face consistency. 40 denoise using Zavy 's excellent ZavyChromaXL v7. DreamshaperXL is really new so this is just for fun. Looking for a tutorial to train your own face using Automatic1111. If the problem still persists I will do the refiner-retraining. The chart above evaluates user preference for SDXL (with and without refinement) over SDXL 0. Second, you'll need the photo style SDXL checkpoint of your preference. At that moment, I was able to just download a zip, type something in webui, and then click generate. Stable Diffusion looks too complicated”. 5 VAE) -> SD 1. Suppose we want a bar-scene from dungeons and dragons, we might prompt for something like. 5 LCM refiner sampler pass. Stable Diffusion is trained on a subset of those images, around 600 million of those, supposedly. There is no such thing as an SD 1. It'll be perfect if it includes upscale too (though I can upscale it in an extra step in the extras tap of automatic1111). Then I do multiple img2img passes with a higher resolution, more VATSIM (Virtual Air Traffic Simulation Network) is the go-to online flight simulation network, where virtual pilots can connect their flight simulators to a shared network and enjoy realistic communication and procedures by VATSIM's trained virtual Air Traffic Controllers. that extension really helps. Use SDXL Refiner with old models. LewdGarlic. Automatic1111 can’t use the refiner correctly. Then play with the refiner steps and strength (30/50 Activate the Face Swapper via the auxiliary switch in the Functions section of the workflow. Technically dreambooth is a also a fine tuning technique. Use a value around 1. I was surprised by how nicely the SDXL Refiner can work even with Dreamshaper as long as you keep the steps really low. 509K subscribers in the StableDiffusion ComfyUI with SDXL (Base+Refiner) + ControlNet XL OpenPose + FaceDefiner (2x) ComfyUI is hard. 74 votes, 16 comments. 6. As per the SD super stage event, the refiner is an optional second pass that can improve some generations. Comparison. I've found very good results doing 15-20 steps with SDXL which produces a somewhat rough image, then 20 steps at 0. Experimental Functions. Karras in general are superb at low step counts (though LMS Karras gets lots of artifacts at high step counts, so never do that). 3) Jul 22, 2023 ยท After Detailer (adetailer) is a Stable Diffusion Automatic11111 web-UI extension that automates inpainting and more. No surprises, Medium is much worse. In my experience, LMS has similar quality to 2M (both Karras and non, compared to their 2M versions), but LMS samplers are more artifact-prone, esp. Try: rear view shot or just rear shot. 1), crowded, alluring eyes, detailed skin, highly detailed, hyperdetailed, intricate, soft lighting, deep focus, photographed on a Canon 5D, 24mm macro lens, F/8 aperture, film still [after]{zoom_enhance mask="face" replacement /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. ๐ท All of the flexibility of Stable Diffusion: SDXL is primed for complex image design workflows that include generation for text or base image, inpainting . a close up of a woman with a butterfly on her head, a photorealistic painting, by Anna Dittmann Flexibility. Medium is using the base workflow on the huggingface Using Refiner -> Base or just CrystalClearXL or other model from the start -> VAEDecode->VAEEncode (SD 1. We'll see about the actual quality, flexibility, prompt adherence and optimization, if/when SD3 comes out fully. 45 denoise it fails to actually refine it. In my opinion the renders of pixart tend to be more interesting and beautiful than the SD3 renders, but they need a second pass with a refiner. Color problems when using face reference. They mostly use python to train Stable Diffusion. Below 0. I see a lot of people complaining about the new hires. Technical details regarding Stable Diffusion samplers, confirmed by Katherine: - DDIM and PLMS are originally the Latent Diffusion repo DDIM was implemented by CompVis group and was default (slightly different update rule than the samplers below, eqn 15 in DDIM paper is the update rule vs solving eqn 14's ODE directly) /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. The main reason why I chose to do this is a selfish one. Install the SDXL auto1111 branch and get both models from stability ai (base and refiner). realvisXL is great and currently probably better than Juggernaut however it is not a Pony model so it can't do what Pony can (but can do a few things that Pony struggles with, like working with controlnet). Step two - upscale: Change the model from the SDXL base to the refiner and process the raw picture in img2img using the Ultimate SD upscale extension with the following settings: (same prompt) Steps: 30, Sampler: DPM++ 2M Karras, CFG scale: 7, Seed: 1799556987, Size: 2304x1792, Model hash: 7440042bbd, Model: sd_xl_refiner_1. I just started learning about Stable Diffusion recently, I downloaded the safe-tensors directly from huggingface for Base and Refiner model, I found…. This simple thing made me a fan of Stable Diffusion. fix with SDXL is broken. Normal Hires. AP Workflow v5. For now, I have to manually copy the right prompts. Key Takeaways. 5 model as your base model, and a second SD1. This series of images is made to see if the color depth in SD3 can be translated in the process of refiner pass. This is not my code, I'm simply posting it. Note: I used a 4x upscaling model which produces a 2048x2048, using a 2x model should get better times, probably with the same effect. SD 1. 0 Base, moved it to img2img, removed the LORA and changed the checkpoint to SDXL 1. And the SDE++ 2M versions are also fast per step. It is suitably sized to become the next standard in text-to-image models. Award. face recognition API. 5 models and LoRA are so fine tuned that while SDXL gives me a much wider range of control, getting the 'perfect' finish seems to only be reliable with 2) Set Refiner Upscale Value and Denoise value. An style can be slightly changed in the refining step, but a concept that doesn't exist in the standard dataset is usually lost or turned into another thing (I. Trained information is represented in an alternate form (using CLIP for text and VAE for image embeddings). Whenever you generate images that have a lot of detail and different topics in them, SD struggles to not mix those details into every "space" it's filling in running through the denoising step. 0 and upscalers. 85, although producing some weird paws on some of the steps. 4 - 0. 0. Even the Comfy workflows aren’t necessarily ideal, but they’re at least closer. Unrealistic body standards having a 24 pack. It is the delicate interplay of shadow and light. Can I use a different CFG value for the refiner in Comfy? I'm currently using Forge, should I switch to Comfy or StableSwarm? 2. It's not, it's just barely better. 5 for final detail refinement seems to give me the ultimate control. 0 of my AP Workflow for ComfyUI. pony is anime and 2d model. , that is more conspicuous than the number of fingers RTX 3060 12GB VRAM, and 32GB system RAM here. ComfyUI with SDXL (Base+Refiner) + ControlNet XL OpenPose + FaceDefiner (2x) ComfyUI is hard. But if you use both together it will make very little differences. AP Workflow 5. The vae that was originally baked into SDXL created visual artifacts when it tried to do its "invisible watermarking". 3), detailed face, /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app So in order to get some answers I'm comparing SDXL1. Legal and PR issues. It detects hands greater than 60x60 pixels in a 512x512 image, fits a mesh model and then generates SDXL vs SDXL Refiner - Img2Img Denoising Plot. Ah also, death to the false emperor, blood for the blood god. 5 and 2. There is an SDXL 0. We need laws that mark images like this as AI generated so we don't get low self-esteem. fix. The negative prompt sounds like frustration to me. If you don't use hires. Here is an example of two images. I just released version 4. 1. Edit: I realized that the workflow loads just fine, but the prompts are sometimes not as expected. This seemed to add more detail all the way up to 0. I fix all my hands in dalle-2. The control Net Softedge is used to preserve the elements and shape, you can also use Lineart) 3) Setup Animate Diff Refiner Hey, a lot of thanks for this! I had a pretty good face upscaling routine going for 1. Reply. Tried a bunch of images, and none of them had the hands detected…. 78. fix @Dr__Macabre. It can do this because it was trained on a lot of images paired with their text captioning with various amount of noise added to the image. Question - Help. 9 safetesnors file. 5 of my wifes face works much better than the ones Ive made with sdxl so I enabled independent prompting(for highresfix and refiner) and use the 1. Thanks for this - newbs coming from A1111 can be overwhelmed by the ComfyUI when trying to locate nodes. If you've seen this post before, you know what to expect. From L to R, this is SDXL Base -- SDXL + Refiner -- Dreamshaper -- Dreamshaper + SDXL Refiner. Every time I use a face reference for my stable diffusion model I get really weird artifacts. These sample images were created locally using Automatic1111's web ui, but you can also achieve similar results by entering prompts one at a time into your distribution/website of choice. 0 and some of the current available custom models on civitai with and without the refiner. Can someone guide me to the best all-in-one workflow that includes base model, refiner model, hi-res fix, and one LORA. but if I run Base model (creating some images with it) without activating that extension or simply forgot to select the Refiner model, and LATER activating it, it gets OOM (out of memory) very much likely when generating images. (as mentioned Used Automatic1111, SDXL 1. 5 model in highresfix with denoise set in the . I know there is the ComfyAnonymous workflow but it's lacking. 5 denoise with SD1. It depends on: what model you are using for the refiner (hint, you don't HAVE to use stabilities refiner model, you can use any model that is the same family as the base generation model - so for example a SD1. Last, I also performed the same test with a resize by scale of 2: SDXL vs SDXL Refiner - 2x Img2Img Denoising Plot. Imagine if you can do your model photoshoot with your new watch, skin care product, or line of overprice handbags in a studio, and seamlessly put the model in the streets of Milan, on the beaches of the Maldives, or wherever else instagram and tiktok says your target demo wants UNet does precisely this, working on different levels of detail as it downscales and upscales. Just like Juggernaut started with Stable Diffusion 1. I want to use Pony as a base model and Juggernaut Lightning as a refiner for more realistic images. Yes. It is not a reasonable approximation, it is the actual data it was trained on. I'll then be wondering why the image was so bad ๐. Edit: RTX 3080 10gb example with a shitty prompt just for demonstration purposes: Without --medvram-sdxl enabled, base SDXL + refiner took 5 mins 6. F222 is a traditional fine tuned model that does not require a special keyword. I will see that when I click on the wrong model and used it instead of the base. SD is a big thing with a lot going on, don't be afraid The truth about hires. I'm using the recommended settings; Sampling Steps: 80-100, Sampler: Euler a, Denoising strength: 0. Switch the timing point from 0. 0 includes the following experimental functions: Free Lunch (v1 and v2) AI researchers have discovered an optimization for Stable Diffusion models that improves the quality of the generated images. But they have different CFG values (because of the lightning). 5 in A1111, but now with SDXL in Comfy I'm struggling to get good results by simply sending an upscaled output to a new pair of base+refiner samplers Code Posted for Hand Refiner. i tried "camera from behind" or "camera shot from behind", i cant really think of other prompts to use, but iv only been able to get like 1 out of 20 images to be from behind with this. The smaller size of this model makes it perfect for running on consumer PCs and laptops as well as enterprise-tier GPUs. It's often not required. You can inpaint with SDXL like you can with any model. text_l & refiner: "(pale skin:1. 5 version, losing most of the XL elements. What they have is a marketable product. It is the curve of rolling hills. To get it back, go to settings --> user interface and add it back. All prompts share the same seed. Reply reply [deleted] High detail RAW color Photo of a strong man, hands in the face, urban city in the background, (full body view:1. If you install locally you can add in your own additions to the main model. Anyway, while i was writing this post, there has been a new update and it now look like this : Here we go. First of all, sorry if this doesn't make sense, i'm french so english isn't my native language and i'm self-taught when it comes to english. The title tells everything. Code for automatically detecting and correcting hands in Stable Diffusion using models of hands, ControlNet, and inpainting. I also automated the split of the diffusion steps between the Base and the Actually the normal XL BASE model is better than the refiner in some points (face for instance) but I think that the refiner can bring some interesting details. Basically it just creates a 512x512 as usual, then upscales it, then feeds it to the refiner. Thanks. Yep! I've tried and refiner degrades (or changes) the results. 9 and Stable Diffusion 1. If you have powerful GPU and 32GB of RAM, plenty of disc space - install ComfyUI - snag the workflow - just an image that looks like this one that was made with Comfy - drop it in the UI - and write your prompt - but the setup is a bit involved - and things don't always go smoothly - you will need the toon model as well - Civitai/HuggingFace I can't get Outpainting to work in Stable Diffusion. 0 model. However, this also means that the beginning might be a bit rough ;) NSFW (Nude for example) is possible, but it's not yet recommended and can be prone to errors. The prompts: (simple background:1. So the website shows all the images SD was trained on and more. A person face changes after ADMIN MOD. 5. We all know SD web UI and ComfyUI - those are great tools for people who want to make a deep dive into details, customize workflows, use advanced extensions, and so on. i dont understand what you need. Use img2img to refine details. I've heard you get better results with full body shots if the source images used for the training were also full body shots, and also keeping the dimension to no more than 512X512 durign generation. Describe the character and add to the end of the prompt: illustration by (Studio ghibli style, Art by Hayao Miyazaki:1. Consistent character faces, designs, outfits, and the like are very difficult for Stable Diffusion, and those are open problems. 8. Steps: (some of the settings I used you can see in the slides) Generate first pass with txt2img The checkboxes (face fix, hires fix) disappeared. 9 (just search in youtube sdxl 0. I haven't played with Dreambooth myself so just going by other people's experience. (viewed from behind:1. next version as it should have the newest diffusers and should be lora compatible for the first time. Not sure if it’s the quality of the image or something but the colors become horrible and the art style becomes much more stylized instead of realistic, even when I try higher resolution images. Is there currently a way to adapt a ComfyUI workflow to avoid the refiner touching any human faces? It's removing details that I want kept there: it makes all faces smooth, de-ages them (I don't want that!) and evens them out, which deletes all of the characters' personalities, age, and uniqueness as a result. For today's tutorial I will be using Stable Diffusion XL (SDXL) with the 0. This simple thing also made my that friend a fan of Stable Diffusion. In any case, we could compare the picture obtained with the correct workflow and the refiner. Uncharacteristically, it's not as tidy as I'd like, mainly due to a challenge I have with passing the checkpoint/model name through reroute nodes. Use 0. So, I'm mostly getting really good results in automatic1111 Yes its human faces only, probably best prompting on your dogs photo using img2img or controlnet. Same with SDXL, you can use any two SDXL models as the base Hopefully Adetailer gets updated soon so you can choose the hands inpaint model instead of inpaint global harmonious. I have tried uninstalling stable diffusion (deleting Taking a good image with a poor face, then cropping into the face at an enlarged resolution of it's own, generating a new face with more detail then using an image editor to layer the new face on the old photo and using img2img again to combine them is a very common and powerful practice. In summary, it's crucial to make valid comparisons when evaluating the SDXL with and without the refiner. fix while using the refiner you will see a huge difference. choose two different styles of models, one as a base, one as a refinement of the model. 5 model as the "refiner"). Interesting, gonna try this tomorrow. 30ish range and it fits her face lora to the image without Very nice. 55 and go from there. The dataset I linked above contains 5 billion images, it's called LAION-5B. It saves you time and is great for quickly fixing common issues like garbled faces. After some testing I think the degradation is more noticeable with concepts than styles. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance. 9 workflow, the one that olivio sarikas video works just fine) just replace the models with 1. If you're using Automatic's GUI there should be an option for full res inpainting so you can mask off the face and generate a new one using a prompt referencing the face and it will generate the face at the full resolution of the image and then scale it down to fit the mask. 3 - 1. I am going to experiment a bit more but if it doesn't work out, I may just use Pixart for the global compositional coherence latent base for SDXL and SD 1. Hi. Stable Diffusion creates images out of pure noise. For prompt use something like face, eye color, hair color, hair style, expression. Setup. Like there's Embeddings, which there are quite a few Me too! realvisXL is awesome at photorealism. 7 in the Denoise for Best results. 5 can be seen in the style of the obvious changes in 0. Simply ran the prompt in txt2img with SDXL 1. Do not use the high res fix section (can select none, 0 steps in the high res section), go to the refiner section instead that will be new with all your other extensions (like control net or whatever other extensions you Using refiner with different settings. E. I downscale my SD pictures before using them in dalle-2, then do img2img again and work with cfg and init strength till they just retouch the dalle2 hands. example of the basic model of photo-realistic style, the model used in the refiner is the anime style. All dreambooth models require a special keyword to condition the image generation but the traditional fine tuning (continue training with a narrow dataset) doesn't. I did run the prompts on the SD3 API again to make sure they haven't changed it, and the results were the same, so it's still good. But I'm not sure what I'm doing wrong, in the controlnet area I find the hand depth model and can use it, I would also like to use it in the adetailer (as described in Git) but can't find or select the depth model (control_v11f1p_sd15_depth) there. Opening the image in stable-diffusion-webui's PNG-info I can see that there are indeed two different sets of prompts in that file and for some reason the wrong one is being chosen. It will just produce distorted, incoherent images. [Cross-Post] I feel the original one is better, high denoise refiner destroys the lighting consistency, especially the hands become flat and even changed the skin color in the third image. 5 refiner node. No Automatic1111 or ComfyUI node as of yet. I have been using automatic1111, don't know much about comfyui. and have to close terminal and restart a1111 again to clear that OOM effect. It adds detail and cleans up artifacts. 2) and used the following negative - Negative prompt: blurry, low quality, worst quality, low resolution, artifacts, oversaturated, text, watermark, logo, signature, out SDXL 1. 9 refiner node. 6 or too many steps and it becomes a more fully SD1. Start with a denoise around . But stable diffusion is faster and I can load the workflow I like. First get a photo of the head in the same direction as the result in mind. 7 and then close to the base model. When SD tries to generate an image that is too different in size and aspect ratio from what it is trained on, you end up getting elongated or multiple features such as two heads, two torsos, and 4 legs. So they put some stuff with the . Here are the solutions: ***Basically, install the refiner extension (sd-webui-refiner). It works to a degree but maybe not enough. 5 secs More than 0. The refiner was trained in tandem with the base, so it will not work without it. SDXL vs DreamshaperXL Alpha, +/- Refiner. i came across the "Refiner extension" in the comments here described as "the correct way to use refiner with SDXL" but i am getting the exact same image between checking it on and off and generating the same image seed a few times as a test. If you aren't using that GUI, the best option is to bring it into GIMP Fooocus-MRE v2. I had the same idea of retraining it with the refiner model and then load the lora for the refiner model with the refiner-trained-lora. Im using automatic1111 and I run the initial prompt with sdxl but the lora I made with sd1. The model doesn’t seem to work for anime images…. If you use ComfyUI you can instead use the Ksampler Regarding the "switching" there's a problem right now with the 1. As you can see the difference is an improvement but the image retains nearly everything. 70 Prompt Comparison: SD3 API vs SD3 Medium. Ideally the refiner should be applied at the generation phase, not the upscaling phase. I created this comfyUI workflow to use the new SDXL Refiner with old models: json here. 01 ~ 1, each increase of 0. I've search but found nothing that seems to use Automatic1111. I'd like to share Fooocus-MRE (MoonRide Edition), my variant of the original Fooocus (developed by lllyasviel ), new UI for SDXL models. Refiner extension not doing anything. 5, we're starting small and I'll take you along the entire journey. Here's how to get the benefits of Pony XL, without the drawbacks of art-style. 0 Base vs Base+refiner comparison using different Samplers. 2. Not too sure how exactly to do all this others that are up to date will know better. 1 sdxl model and 1 sd1. 7 in the Refiner Upscale to give a little room in the image to add details. 0 base model and HiresFix x2. The soft inpainting feature is also handy, it tends to blend the seams very well on the inpainted area. 2), cottage. Forcing Lora weights higher breaks the ability for generalising pose, costume, colors, settings etc. Stable Diffusion 3 Medium is Stability AI’s most advanced text-to-image open model yet, comprising two billion parameters. Then install the SDXL Demo extension . Structured Stable Diffusion courses. In this post, you will learn how it works, how to use it, and some common use cases. 9 vae, along with the refiner model. You'll need two checkpoints "Real Pony" checkpoint (there are three, you want the one with the most upvotes on civit to start): this will be the primary checkpoint. You just can't change the conditioning mask strength like you can with a proper inpainting model, but most people don't even know what that is. Hires fix is still there, you just need to click to expand but face restore has indeed been removed from the main page. Generate your images through automatic1111 as always and then, go to the SDXL Demo extension tab, turn on 'Refine' checkbox and drag your image onto the square. Workflow Overview: txt2Img API. xoyjpnnmfxvncthtdhfg